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Image Generation ​

The Image Studio provides a local environment to generate high-resolution images. You can generate images through Stable Diffusion WebUI Forge or ComfyUI.

The studio supports FLUX.1, SDXL, and Stable Diffusion 1.5 base models.

Supported Model Families ​

Select the model family that matches your GPU memory and artistic goals:

Model FamilyNative ResolutionMinimum VRAMRecommended Use Case
FLUX.1 (dev / schnell)1024 × 102412 GBBest text rendering, complex prompts, and realistic anatomy.
SDXL (Stable Diffusion XL)1024 × 10248 GBPhotorealism, cinematic lighting, and custom community styles.
Stable Diffusion 1.5512 × 5124 GBHigh-speed rendering and lightweight game assets.

NOTE

You can download base models and checkpoints directly through the Model Management tab from Hugging Face or CivitAI.

Generation Procedure ​

Follow these steps to generate an image:

Step 1: Check Backend Engine Status ​

  1. Open the Workflows tab in the top navigation bar.
  2. Click 🎨 Images in the modality selector.
  3. Check the engine status pill in the header.
  4. Verify that the status shows Online.
  5. Click Toggle Forge Engine if the status shows Offline.

Step 2: Choose a Preset or Starter Style ​

Select a preset from the Studio Preset Bar to configure optimal resolutions automatically:

  • Standard Square (1024 × 1024): Default setting for general art and character portraits.
  • Landscape Wallpaper (1344 × 768): Wide aspect ratio for environments and desktop backgrounds.
  • Portrait Photo (768 × 1152): Vertical aspect ratio for full-body human figures and posters.
  • Classic SD 1.5 (512 × 512): Lightweight resolution for legacy models and rapid drafting.

TIP

Click a starter prompt chip above the prompt box. The chip populates proven prompts and selects the matching preset instantly.

Step 3: Configure Generation Parameters ​

Set the parameters in the configuration panel:

  1. Workflow Model: Type or select your checkpoint name (for example, SDXL Base or flux1-dev).
  2. Positive Prompt: Describe all subjects, styles, lighting, and camera angles that you want in the image.
  3. Negative Prompt: List unwanted elements such as blur, distortion, or watermark artifacts.
  4. Resolution (Width × Height): Set the image dimensions. Keep dimensions aligned to multiples of 64 or 16.
  5. Sampling Steps:
    • Use 4 to 8 steps for FLUX.1-schnell.
    • Use 20 to 30 steps for SDXL and SD 1.5.
    • Use 25 to 50 steps for FLUX.1-dev.
  6. CFG Scale (Guidance Scale):
    • Set 1.0 for FLUX models.
    • Set 5.0 to 7.0 for SDXL models.
    • Set 7.0 to 8.5 for SD 1.5 models.
  7. Sampler and Scheduler: Select an approved sampler algorithm:
    • Euler or Euler a: Fast, general-purpose convergence.
    • DPM++ 2M Karras: High detail for photorealistic textures.
    • UniPC: High-speed generation with fewer steps.
  8. Seed: Enter a numeric value between 1 and 999999999. Enter -1 or leave random to generate new variations.

Step 4: Launch Generation ​

  1. Click 🎨 Launch Forge Image Studio (Port 7860) or Generate Image.
  2. Observe the 4-stage progress tracker during diffusion execution.
  3. Inspect the completed image in the output gallery.

IMPORTANT

The application automatically saves rendered images to outputs/images/ with embedded prompt and parameter metadata.

Applying LoRAs (Low-Rank Adaptations) ​

LoRAs are compact model files (typically 50 MB to 200 MB). A LoRA modifies an existing checkpoint into a specific character, clothing style, or visual aesthetic.

How to Download LoRAs ​

  1. Open the Stable Diffusion (CivitAI) tab.
  2. Set the Type dropdown filter to LoRA.
  3. Sort results by Highest Rated or Most Downloaded.
  4. Type your desired aesthetic into the search bar (for example, Pixel Art XL or Detail Tweaker).
  5. Click ⬇ Download to Forge. The system stores the file in your models directory.
Style CategorySearch TermRecommended Base ModelEffect
Pixel ArtPixel Art XLSDXLGenerates authentic 16-bit sprites and backgrounds.
Retro 3DPS1 GraphicsSDXL / SD 1.5Recreates jagged low-polygon PlayStation 1 aesthetics.
Anime & CartoonCel ShadedSDXL / Pony V6Applies clean outlines and flat anime color fills.
Micro-DetailDetail Tweaker XLSDXL / FLUXSharpens skin pores, fabric weaves, and rim lighting.

Method 1: Using LoRAs in Prompts (Forge & WebUI) ​

Add the trigger tag directly into your positive prompt box:

text
<lora:pixel_art_xl:0.8> a cyberpunk street market at night, 16-bit retro game asset
  • Specify the filename of the LoRA between the colons.
  • Add a weight value at the end (typically 0.6 to 1.0).
  • Lower the weight if the style distorts facial anatomy or colors.

Method 2: Using LoRAs in ComfyUI ​

When using ComfyUI as your studio backend:

  1. Add a Load LoRA node to your ComfyUI canvas.
  2. Connect the MODEL output from your Checkpoint Loader into the model input of the LoRA node.
  3. Connect the CLIP output into the clip input of the LoRA node.
  4. Select your downloaded LoRA filename inside the node dropdown menu.
  5. Set strength_model and strength_clip to 0.8.
  6. Connect the modified outputs into your positive prompt and KSampler nodes.

VRAM Management and Hardware Fit ​

High-resolution diffusion models require dedicated GPU VRAM. Review these memory guidelines:

  • Automated LLM Offload: The VRAM Orchestrator automatically unloads Ollama models before diffusion starts.
  • FLUX Models: FLUX.1-dev requires at least 12 GB VRAM in FP8 or NF4 quantization.
  • Resolution Limits: Do not set resolutions above 1024 × 1024 on GPUs with 8 GB VRAM. Use high-resolution fix or latent upscalers instead.

WARNING

Generating at resolutions above 1536 × 1536 without tiling causes Out-Of-Memory errors on 16 GB GPUs. Use the standard presets to maintain safe VRAM bounds.

Released under the MIT License.

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